TL;DR:
torch.compile
. The following are benchmarks on one RTX 4090 GPU:Setup | one full training loop (ms) |
---|---|
This repo | 142.44 |
PyTorch | 66.73 |
PyTorch with torch.compile | 59.20 |
To train a diffusion model in CUDA with some sample images from ImageNet 64x64, run the following:
gunzip data/elephant_train.bin.gz # prepare the data python train_unet.py --init_model_only True # need to initialize model weights via python make train_unet ./train_unet
To train the model with your own data, you need to create a .bin
file with your data first:
python prepare_data.py --data_dir YOUR_DATA_DIR --output_name YOUR_BINARY_DATA_FILENAME.bin # now run training, assuming you have already initialized the model as above ./train_unet --data_file YOUR_BINARY_DATA_FILENAME.bin
The PyTorch training code is essentially taken from the guided-diffusion repo. To run PyTorch training, do:
python train_unet.py --data_dir YOUR_DATA_DIR # use --compile 0 if you don't want to call torch.compile() on the model
The CUDA training loop will save model weights in .bin
files. To generate new images with model weights saved in either .bin
or .pt
files, run:
python generate.py --model_filename YOUR_MODEL_WEIGHTS_FILENAME
Inspired by Andrej Karpathy's llm.c, I built a UNet from scratch in C/CUDA. The goal of the project is to learn the concepts in llm.c, and to reach for PyTorch's performance with our CUDA implementation. I chose the UNet because it is a key architecture for diffusion models, and I will do some simple diffusion model training with it.
Diffusion model training is quite sophisticated nowadays. Since this project is focused on learning CUDA as opposed to building the best diffusion model, I prioritized simplicity over performance, and followed the implementation from the paper Diffusion Models Beat GANs on Image Synthesis. Currently the UNet only supports unconditioned diffusion training. I also did not reproduce all the model configurations from the paper; the details of the differences will be explained in the section on the architecture.
Here are some images generated with our CUDA implementation. The model is trained on elephant images from ImageNet 64x64 without class-conditioning. The model is highly over fitting the training set right now, but at least this confirms training is working.
<p align="center"> <img src="assets/cuda_sample1.jpg" alt="Sample 1" width="20%" /> <img src="assets/cuda_sample2.jpg" alt="Sample 2" width="20%" /> <img src="assets/cuda_sample3.jpg" alt="Sample 3" width="20%" /> </p>The Github repository is organized as follows:
dev/
directory contains all different kernels and tests written during development.
.cu
file (e.g. groupnorm.cu
), which contains different CUDA kernels for a layer, and a .py
file (e.g. groupnorm.py
), which contains an identical Pytorch implementation for the same layer.train_unet.cu
is a single file with the full diffusion model training code (~ 5000 lines). We take the best kernels from dev/
and copy them here. The file also contains things like the data loader and AdamW.For a tutorial on how to write the forward and backward pass of different layers, I recommend Andrej's layernorm tutorial.
The rest of these notes are organized as follows. The next section will cover some background, both on diffusion models and on the UNet architecture we use. Then the later sections document successive iterations on the model where I benchmark kernels and try to speed things up. It turns out that most of a UNet's running time is spent doing 3x3
image convolutions, so that is where most of the work went into and where these notes focus on.
Our goal is to train a diffusion model with a UNet. Let me give a short summary of how diffusion models work. A good mathematical description can be found in Appendix B of the paper Diffusion Models Beat GANs on Image Synthesis; a good hands-on tutorial can be found at Chenyang Yuan's blog. We start with a target distribution $\pi(x)$ on $\mathbb{R}^d$ that we want to sample from. In our case, the space will be RGB images with $C = 3$ channels, height and weight $H = W = 64$, and $d = C \times H \times W$, and the target distribution will be elephant images. The key idea is to set up a stochastic process $(X_t)_{t \ge 0}$ with the following three properties:
These properties together enable us to draw samples from the target $\pi$ as follows:
So now we need a stochastic process that satisfies these properties, and a way to learn the conditional distributions in property 3. The stochastic process $(X_t)_{t \ge 0}$ will look like so: $X_0$ is drawn from $\pi$, and $X_t$ is distributed as follows:
$$ X_t = \sqrt{\alpha_t} \cdot X_0 + \sqrt{1 - \alpha_t}\cdot \epsilon, $$
where $\epsilon$ is a standard Gaussian in $\mathbb{R}^d$, and $\alpha_t$ is a non-increasing function of $t$ that we will choose, with the properties that $\alpha_0 = 1$ and $\alpha_t \to 0$ as $t \to \infty$. We see that when $t$ is large, $X_t \approx \epsilon$, which satisfies property 2. Note that the equation above is only specifying the marginal distribution of $X_t$, so the conditional distribution $\pi(X_{t-1} \mid X_t)$ may not be deterministic (when the conditional is deterministic, we have the DDIM models).
To sample from the conditional distribution $\pi(X_{t-1} \mid X_t)$, we will train a model $\epsilon_\theta(X_t, t)$ that takes $X_t$ and $t$ as input and minimizes the following objective:
$$ L = \mathbb{E}[\lVert \epsilon - \epsilon_\theta(X_t, t) \rVert^2]. $$
Here the expectation is taken over $\epsilon$, $t$ and $X_t$, where $t$ uniformly sampled from the range $[0, T]$, $\epsilon$ is sampled from the standard Gaussian, $X_0$ is sampled from $\pi$ (i.e. one of our training data), and $X_t$ is then constructed from $X_0$, $t$, and $\epsilon$ using the identity above. Conceptually the model $\epsilon_\theta$ takes in the noisy input $X_t$, and tries to learn the noise component $\epsilon$ within the input. With this model, it is then fairly easy to do the conditional sampling from $\pi(X_{t - 1} \mid X_t)$; the details can be found in Appendix B of Diffusion Models Beat GANs on Image Synthesis.
Our loss function dictates that we want a model which takes an input of shape (B, C, H, W)
, where B
is the batch dimension, and returns an output of the same shape. The UNet is a sample efficient architecture designed specifically for such scenarios. The UNet we use is a basic version taken from the paper Diffusion Models Beat GANs on Image Synthesis, and it looks like this:
Specifically, we use the residual blocks from BigGAN, which look like so (from Figure 15 of Large Scale GAN Training for High Fidelity Natural Image Synthesis ):
<p align="center"> <img src="assets/resblock.png" width="30%" /> </p>A few more notes on model details:
train_unet.py
. Our model exactly matches the official implementation with the following model configurations:--attention_resolutions 16,8 --class_cond False --diffusion_steps 1000 --dropout 0.0 --image_size 64 --learn_sigma False --noise_schedule linear --num_channels 64 --num_head_channels 32 --num_res_blocks 2 --resblock_updown False --use_new_attention_order True --use_fp16 False --use_scale_shift_norm False
In the first version I wanted to get something working quickly, so I copied or adapted the kernels from llm.c. This approach took care of some kernels: the linear layer could be reused from llm.c without adaption; the groupnorm layer is different from the layernorm in llm.c, but we only needed to change the axis we reduce over and then we had a working kernel.
The self-attention layer was trickier. At first glance the adaptation seems straightforward: the attention layer functions identically for transformers and image models, and the only difference is that instead of the inputs having shape (B, T, C)
, they now have (B, C, H, W)
. So we can reuse the transformer attention kernels by first transposing the inputs to shape (B, H * W, C)
, then calling the kernels with T = H * W
, then transposing the output back to shape (B, C, H, W)
.
This turns out to be highly inefficient, because for each transpose we need to move a block of size B * C * H * W
in and out of GPU global memory, and as we will see later such steps should be avoided. So the attention kernels will be an obvious place for future improvements.
Several kernels did not exist in llm.c, but they are needed for the UNet. They are:
3x3
and 1x1
convolutions.The up and down sample kernels (nearest interpolation and average pooling respectively) are easy: there is barely any computation, and we easily parallelize them by assigning one pixel to each GPU thread.
So we are left with the convolution kernels. I wanted to get something working quickly, but I also didn't want it to be too slow, so my plan was to convert all the convolutions to matrix multiplications, and then use cuBLAS, which should be fast.
This plan is quite natural for the 1x1
convolution: for inputs of shape (B, C_in, H, W)
and weights of shape (C_out, C_in)
, the forward pass for a 1x1
convolution is essentially a matrix multiplication in the C_in
dimension of the input with the weights. So 1x1
convolutions are done with the following steps:
(B, C_in, H, W)
to (B * H * W, C_in)
,(B * H * W, C_out)
, then add the bias,(B, C_out, H, W)
.Notice again that this approach needs two transposes of the entire input array, which are expensive. In iteration 2 we will write a custom kernel that avoids these transposes.
For the 3x3
convolutions, things are trickier. Let's focus on the forward pass, where the shapes of the relevant parameters are as follows:
(B, C_in, H, W)
,(C_out, C_in, 3, 3)
,(B, C_out, H, W)
.Since the plan is to cast the convolution into a matmul, it seems natural to transpose the input and output to shapes (B * H * W, C_in)
and (B * H * W, C_out)
respectively. For the weight tensor, we can think of it as consisting of 9 different weight matrices, all of shape (C_out, C_in)
, where each one corresponds to one of the 9 filters in the 3x3
convolution. Let's introduce some notation: let the transposed input be $X \in \mathbb{R}^{(B\cdot H \cdot W) \times C_\text{in}}$, the transposed output be $Y \in \mathbb{R}^{(B\cdot H\cdot W) \times C_\text{out}}$, and let the weight tensor be $W\in \mathbb{R}^{C_\text{out} \times C_\text{in} \times 9}$, where $W = (W_0, \dots, W_8)$, and each $W_i \in \mathbb{R}^{C_\text{out} \times C_\text{in}}$ is the weight matrix for filter $i \in {0, \dots, 8}$.
The convolution of a single pixel works by multiplying the pixels with the 9 filters and summing the values. So roughly speaking, the convolution for the entire batch looks something like this:
(B * H * W, C_out)
.So we have turned the 3x3
convolution into matrix multiplications. Except this is not quite right, because not every pixel is multiplied with every filter. For instance, if we
一键生成PPT和Word,让学习生活更轻松
讯飞智文是一个利用 AI 技术的项目,能够帮助用户生成 PPT 以及各类文档。无论是商业领域的市场分析报告、年度目标制定,还是学生群体的职业生涯规划、实习避坑指南,亦或是活动策划、旅游攻略等内容,它都能提供支持,帮助用户精准表达,轻松呈现各种信息。
深度推理能力全新升级,全面对标OpenAI o1
科大讯飞的星火大模型,支持语言理解、知识问答和文本创作等多功能,适用于多种文件和业务场景,提升办公和日常生活的效率。讯飞星火是一个提供丰富智能服务的平台,涵盖科技资讯、图像创作、写作辅助、编程解答、科研文献解读等功能,能为不同需求的用户提供便捷高效的帮助,助力用户轻松获取信息、解决问题,满足多样化使用场景。
一种基于大语言模型的高效单流解耦语音令牌文本到语音合成模型
Spark-TTS 是一个基于 PyTorch 的开源文本到语音合成项目,由多个知名机构联合参与。该项目提供了高效的 LLM(大语言模型)驱动的语音合成方案,支持语音克隆和语音创建功能,可通过命令行界面(CLI)和 Web UI 两种方式使用。用户可以根据需求调整语音的性别、音高、速度等参数,生成高质量的语音。该项目适用于多种场景,如有声读物制作、智能语音助手开发等。
字节跳动发布的AI编程神器IDE
Trae是一种自适应的集成开发环境(IDE),通过自动化和多元协作改变开发流程。利用Trae,团队能够更快速、精确地编写和部署代码,从而提高编程效率和项目交付速度。Trae具备上下文感知和代码自动完成功能,是提升开发效率的理想工具。
AI助力,做PPT更简单!
咔片是一款轻量化在线演示设计工具,借助 AI 技术,实现从内容生成到智能设计的一站式 PPT 制作服务。支持多种文档格式导入生成 PPT,提供海量模板、智能美化、素材替换等功能,适用于销售、教师、学生等各类人群,能高效制作出高品质 PPT,满足不同场景演示需求。
选题、配图、成文,一站式创作,让内容运营更高效
讯飞绘文,一个AI集成平台,支持写作、选题、配图、排版和发布。高效生成适用于各类媒体的定制内容,加速品牌传播,提升内容营销效果。
专业的AI公文写作平台,公文写作神器
AI 材料星,专业的 AI 公文写作辅助平台,为体制内工作人员提供高效的公文写作解决方案。拥有海量公文文库、9 大核心 AI 功能,支持 30 + 文稿类型生成,助力快速完成领导讲话、工作总结、述职报告等材料,提升办公效率,是体制打工人的得力写作神器。
OpenAI Agents SDK,助力开发者便捷使用 OpenAI 相关功能。
openai-agents-python 是 OpenAI 推出的一款强大 Python SDK,它为开发者提供了与 OpenAI 模型交互的高效工具,支持工具调用、结果处理、追踪等功能,涵盖多种应用场景,如研究助手、财务研究等,能显著提升开发效率,让开发者更轻松地利用 OpenAI 的技术优势。
高分辨率纹理 3D 资产生成
Hunyuan3D-2 是腾讯开发的用于 3D 资产生成的强大工具,支持从文本描述、单张图片或多视角图片生成 3D 模型,具备快速形状生成能力,可生成带纹理的高质量 3D 模型,适用于多个领域,为 3D 创作提供了高效解决方案。
一个具备存储、管理和客户端操作等多种功能的分布式文件系统相关项目。
3FS 是一个功能强大的分布式文件系统项目,涵盖了存储引擎、元数据管理、客户端工具等多个模块。它支持多种文件操作,如创建文件和目录、设置布局等,同时具备高效的事件循环、节点选择和协程池管理等特性。适用于需要大规模数据存储和管理的场景,能够提高系统的性能和可靠性,是分布式存储领域的优质解决方案。
最新AI工具、AI资讯
独家AI资源、AI项目落地
微信扫一扫关注公众号